Not looking to change backgrounds or anything particularly, just the character (realistically).

Is any way to save whole workflow in a1111 or Forge with control net images?

I tried CogVideoX with starting frame I2V and it was great. I'm not sure if you can hack start and end frames with it yet. I know DynamiCrafter Interpolation is there, but its U-Net based and I'm looking for DiT based models.

There is no shortage of photorealistic checkpoints, but when i have to pick an "artistic" alrounder -no anime- for sdxl seems like a more difficult choice. it's still Juggernaut the best choice? zavychroma? albedobase?

I'd like to read your suggestions.

Been playing around in Forge for a few weeks now and I finally decided to jump into text2video to see what my technically illiterate self could do. Unfortunately, while most extensions seem cross-compatible, text2video gives me a bunch of errors and won't generate a tab for itself.

Is there an alternative I need to grab for Forge, or should I just install A1111 on the side for that purpose?

Edit:

So apparently this is actually a general error, since I'm getting the same error on my fresh installation of A1111:

Error loading script: api_t2v.py

Traceback (most recent call last):

File "B:\AI\Stability Matrix\Data\Packages\stable-diffusion-webui-forge\modules\scripts.py", line 525, in load_scripts

script_module = script_loading.load_module(scriptfile.path)

File "B:\AI\Stability Matrix\Data\Packages\stable-diffusion-webui-forge\modules\script_loading.py", line 13, in load_module

module_spec.loader.exec_module(module)

File "<frozen importlib._bootstrap_external>", line 883, in exec_module

File "<frozen importlib._bootstrap>", line 241, in _call_with_frames_removed

File "B:\AI\Stability Matrix\Data\Packages\stable-diffusion-webui-forge\extensions\sd-webui-text2video\scripts\api_t2v.py", line 39, in <module>

from t2v_helpers.args import T2VArgs_sanity_check, T2VArgs, T2VOutputArgs

File "B:\AI\Stability Matrix\Data\Packages\stable-diffusion-webui-forge/extensions/sd-webui-text2video/scripts\t2v_helpers\args.py", line 7, in <module>

from samplers.samplers_common import available_samplers

File "B:\AI\Stability Matrix\Data\Packages\stable-diffusion-webui-forge/extensions/sd-webui-text2video/scripts\samplers\samplers_common.py", line 2, in <module>

from samplers.ddim.sampler import DDIMSampler

File "B:\AI\Stability Matrix\Data\Packages\stable-diffusion-webui-forge/extensions/sd-webui-text2video/scripts\samplers\ddim\sampler.py", line 7, in <module>

from ldm.modules.diffusionmodules.util import make_ddim_sampling_parameters, make_ddim_timesteps, noise_like, extract_into_tensor

ModuleNotFoundError: No module named 'ldm'

Error loading script: text2vid.py

Traceback (most recent call last):

File "B:\AI\Stability Matrix\Data\Packages\stable-diffusion-webui-forge\modules\scripts.py", line 525, in load_scripts

script_module = script_loading.load_module(scriptfile.path)

File "B:\AI\Stability Matrix\Data\Packages\stable-diffusion-webui-forge\modules\script_loading.py", line 13, in load_module

module_spec.loader.exec_module(module)

File "<frozen importlib._bootstrap_external>", line 883, in exec_module

File "<frozen importlib._bootstrap>", line 241, in _call_with_frames_removed

File "B:\AI\Stability Matrix\Data\Packages\stable-diffusion-webui-forge\extensions\sd-webui-text2video\scripts\text2vid.py", line 24, in <module>

from t2v_helpers.render import run

File "B:\AI\Stability Matrix\Data\Packages\stable-diffusion-webui-forge/extensions/sd-webui-text2video/scripts\t2v_helpers\render.py", line 5, in <module>

from modelscope.process_modelscope import process_modelscope

ModuleNotFoundError: No module named 'modelscope.process_modelscope'

I've been at this for way too long. I'm trying to get it to train a LoRA for Flux - I had several issues with getting it to actually install the Flux branch of Kohya_SS, but it eventually installed supposedly fine. The issue is that it doesn't start training because it thinks there is something wrong with the VAE: https://huggingface.co/black-forest-labs/FLUX.1-dev/blob/main/vae/diffusion_pytorch_model.safetensors

I tried with both the sd3-flux.1 branch and the sd3-sd3.5-flux branch, and they both produce this log:

INFO     Loading state dict from                                                 flux_utils.py:149
                                   
C:/ComfyUI/ComfyUI_windows_portable/ComfyUI/models/vae/model.FLUX.1dev.safetensors

INFO     Loaded AE:                                                              flux_utils.py:152
         _IncompatibleKeys(missing_keys= [bunch of encoders/decoders])

And

Traceback (most recent call last):
  File "C:\Kohya_SS\kohya_ss\sd-scripts\flux_train.py", line 998, in <module>
    train(args)
  File "C:\Kohya_SS\kohya_ss\sd-scripts\flux_train.py", line 189, in train
    ae.to(accelerator.device, dtype=weight_dtype)
  File "C:\Kohya_SS\kohya_ss\venv\lib\site-packages\torch\nn\modules\module.py", line 1340, in to
    return self._apply(convert)
  File "C:\Kohya_SS\kohya_ss\venv\lib\site-packages\torch\nn\modules\module.py", line 900, in _apply
    module._apply(fn)
  File "C:\Kohya_SS\kohya_ss\venv\lib\site-packages\torch\nn\modules\module.py", line 900, in _apply
    module._apply(fn)
  File "C:\Kohya_SS\kohya_ss\venv\lib\site-packages\torch\nn\modules\module.py", line 900, in _apply
    module._apply(fn)
  [Previous line repeated 3 more times]
  File "C:\Kohya_SS\kohya_ss\venv\lib\site-packages\torch\nn\modules\module.py", line 927, in _apply
    param_applied = fn(param)
  File "C:\Kohya_SS\kohya_ss\venv\lib\site-packages\torch\nn\modules\module.py", line 1333, in convert
    raise NotImplementedError(
NotImplementedError: Cannot copy out of meta tensor; no data! Please use torch.nn.Module.to_empty() instead of torch.nn.Module.to() when moving module from meta to a different device.
Traceback (most recent call last):
  File "C:\Users\user\AppData\Local\Programs\Python\Python310\lib\runpy.py", line 196, in _run_module_as_main
    return _run_code(code, main_globals, None,
  File "C:\Users\user\AppData\Local\Programs\Python\Python310\lib\runpy.py", line 86, in _run_code
    exec(code, run_globals)
  File "C:\Kohya_SS\kohya_ss\venv\Scripts\accelerate.EXE\__main__.py", line 7, in <module>
    sys.exit(main())
  File "C:\Kohya_SS\kohya_ss\venv\lib\site-packages\accelerate\commands\accelerate_cli.py", line 48, in main
    args.func(args)
  File "C:\Kohya_SS\kohya_ss\venv\lib\site-packages\accelerate\commands\launch.py", line 1106, in launch_command
    simple_launcher(args)
  File "C:\Kohya_SS\kohya_ss\venv\lib\site-packages\accelerate\commands\launch.py", line 704, in simple_launcher
    raise subprocess.CalledProcessError(returncode=process.returncode, cmd=cmd)
subprocess.CalledProcessError: Command '['C:\\Kohya_SS\\kohya_ss\\venv\\Scripts\\python.exe', 'C:/Kohya_SS/kohya_ss/sd-scripts/flux_train.py', '--config_file', 'C:/Kohya_SS/kohya_ss/outputs\\model/config_dreambooth-20241226-100738.toml']' returned non-zero exit status 1.
10:07:46-798363 INFO     Training has ended.

Is it something with the VAE version? Is it something with the Python installation? It seems no one has raised this specific bug on the project.

Thank you in advance.

hello everyone, I need some advice on how to generate photorealistic images of hand grasping objects given the object image. Given only an white-bkgd image of the object like the first image below, how to generate a good one like the second below?

Should I use img2img, img inpainting or image editing method?

Thank you for your advice!

input: object img

output: hand object interaction

Kinda new with stable diffusion, but reading on some subs, I noticed that people are generating images larger than the max resolution of models. I am doing image2image generation i.e. input is prompt + image. Right now this does not work with large images, like lets say I wanted generate images of a website banner for example. I do have a subscription with stability (dont have my own compute or means to it), how would I go about doing this?

Hello everyone!

I'm reaching out to the community for some guidance on a technique I've been trying to master using SDXL or FLUX. My goal is to take an existing concept image and replace specific areas with different images, while keeping the overall composition intact.

Example:

For instance, I have a concept photo featuring a pair of jeans, and I want to isolate the jeans area and replace it with chinos, maintaining the same pose and background.

My Experience:

I've experimented with various methods, including SD1.5 DreamBooth, SDXL DreamBooth LoRA, IPAdapter, and others, but I haven't been able to achieve the results I want. I would like to train with 3 to 4 target images so that I can generate this object using just a token prompt, while also capturing details like buttons and other intricate features.

Questions:

What are the best practices for masking and replacing specific areas in an image using SDXL or FLUX?

Are there any specific prompts or settings that have worked well for you in achieving seamless image replacements?

How can I effectively train the model with a few target images to ensure I get the desired output?

I’d really appreciate any tips, techniques, or resources you could share! Thank you in advance for your help!

Feel free to modify anything to better!

https://preview.redd.it/eczc843l459e1.png?width=435&format=png&auto=webp&s=233538fcb4f5187d2453ea68b71917e068d4f9da

Hi everyone,

I recently installed the Krita AI Diffusion plugin using the guidelines provided here. While the plugin is working to some extent, I've noticed that several options in the AI Image Generation dropdown menu are missing. Features like "Expand," "Add Content," "Remove," "Replace," and "Fill" aren't showing up.

Has anyone else experienced this issue? Could it be related to the installation process, dependencies, or perhaps my version of Krita? I'd appreciate any advice or troubleshooting tips to get those missing features to appear.

Thanks in advance!

I've got this image of Santa Claus and alien elves, but when I used LTX-Video with the prompt "Santa Claus walking and nodding towards alien elves," the output was pretty bad. Here's the image I'm working with: Image.

Can anyone recommend some other tools or methods to create a decent motion video from this static image? I'm looking for something that can handle this unusual scene well.

I've heard or read somewhere that comfy can only utilize Nvidia cards. This obviously limits selection quite heavily, especially with cost in mind. Is this information accurate?

I'm trying to understand what is the issue. I have the "Image folder (containing training images subfolders)" set to "C:/kohya_ss/training_images" and the "Training images (directory containing the training images)" set to "C:/kohya_ss/training_images/name". It keeps loading the upper folder successfully but then say it skips the actual image-containing folder because it doesn't contain an underscore.

I don't know why that would matter but I tried adding an underscore to the folder name, naming the files inside of it instead of "1.png", "2.png" etc. "name_1.png" etc, it doesn't work.

That error is not marked red, but after loading the other stuff it has a red error "No data found. Please verify arguments (train_data_dir must be the parent of folders with images)" and the training is aborted.

How do I fix this? Thank you in advance.

I decided to try the new SD 3.5 medium, coming from the SDXL models, I think the SD 3.5 medium has a great potential, much better compared to the base SDXL model, even comparable to fine-tuned SDXL models.

Since I don´t have a beast GPU, just my personal laptop, takes up to 3 minutes to generate with Flux models, but SD 3.5 medium is a nice spot between SDXL and FLUX.

I combined the turbo and 3 small LORAs and got good results with 10 steps:

WORKFLOW: https://civitai.com/posts/10757286

### 1

Dark Maccabre Art, Gothic Horror, Creepy Demonic Witch. Faceless. Hooded. Long Purple Hair. Veil created from thick fog. she is holding a sphere of mesmerzing mana in her hands. glowing particles. ultrarealistic and detailed. 8K

### 2

a striking and surreal scene that combines elements of both the natural world and fantasy. Dominating the composition is a massive, reptilian eye, filling almost the entire frame. The eye is highly detailed, with a slit-like pupil that suggests it belongs to a large, powerful creature, perhaps a dragon or another mythical being. The texture around the eye is rugged and scaly, giving the impression of ancient, weathered skin. In the lower portion of the image, a solitary human figure stands before the eye, dressed in a flowing black robe. The figure is tiny in comparison to the colossal eye, emphasizing the vast difference in scale and power between the two. The person stands on a surface that appears to be water or mist, which reflects the eerie, otherworldly light that surrounds the scene. The atmosphere is misty and dreamlike, adding to the sense of mystery and awe. Overall, the image is both dramatic and thought-provoking, blending cultural elements with a fantastical imagination to create a visually captivating scene.

### 3

A breathtaking sunset panorama painting in style of Van Gogh and Nicholas Roerich of a tropical beach on Ganymede, Jupiter in the night sky, cerulean and maroon palette, impressionism,

### 4

A Closeup Portrait of an DARK Arab girl, extreme Closeup of her Face - shrouded in mystery. She wears a, tattered high Arabic patterns scarf in a mesmerizing blend of vibrant colors, including neon pink, blue, green, and purple, which create an otherworldly, glowing effect. The fabric seems to blend seamlessly with the natural environment, as if it's a part of the sky. Hyperdetailed badass Closeup, hyperdetailed, deadly Gaze, mouth obscured by the coats high collar

### 5

a dark fantasy portrait of a powerful frozen necromancer emerging from swirling froze and embers. The necromancer should have dark energy of ice, cracked ice skin, glowing blue sockets in scull under hood. Its expression should be menacing and powerful. The background should be filled with dark, swirling smoke interwoven with bright blue embers. Use dramatic lighting to highlight the necromancer's features and create a sense of depth. The overall mood should be dark, ominous, and terrifying. The style should be reminiscent of dark fantasy illustrations with a high level of detail and realism. Aim for a cinematic, impactful composition with a shallow depth of field, focusing on the necromancer's scull. The color palette should be limited to dark blues of scull and embers.

### 6

the lady of the golden hour by Russ Mills

### 7

8k, UHD, best quality, highly detailed, cinematic, photographic, a female space soldier wearing an orange and white space suit exploring a river in a dark mossy canyon on another planet, full body photo away from camera, helmet, gold tinted face shield, (glowing fireflies), (dark atmosphere), haze, halation, bloom, dramatic atmosphere, sci-fi movie still, (jungle), (moss)

### 8

Oil painting by Montague Dawson titled \"The Stately Ship.\" Depicts a full-rigged ship sailing on a turbulent sea. Ship centered in composition, angled slightly to the right, showcasing detailed sails and rigging catching the wind. Blue waves with whitecaps occupy the foreground, suggesting movement and depth. Horizon line low, allowing expansive sky with soft clouds. Lighting suggests early morning or afternoon with soft shadows. Art style falls under marine art, capturing dynamic realism and meticulous attention to nautical detail. Signature in the lower left.

### 9

a highly detailed realistic CGI rendered image in a fantasy style, depicting a whimsical winter forest scene. At the center of the image is an owl with large, expressive brown eyes, sitting on a moss-covered rock. The owl is wearing a green knitted beanie hat, adding a touch of charm and personality. Its feathers are a mix of white and brown, blending seamlessly into the snowy environment. Surrounding the owl are various elements that enhance the magical atmosphere. To the left of the owl, a large, bright orange mushroom with a white cap covered in snow stands tall on a tree stump. The mushroom emits a soft, warm light, contrasting with the cool, wintry tones of the scene. In the background, the forest is filled with tall, snow-covered trees, their branches bare and twisted, creating a mysterious and enchanting backdrop. The ground is blanketed with fresh snow, and the forest floor is dotted with glowing, luminescent mushrooms, adding a mystical touch. The lighting in the image is soft and diffused, with a gentle glow from the mushrooms and the mushroom cap, creating a serene and magical winter wonderland. The overall mood is peaceful and enchanting, inviting viewers into a fantastical world.

### 10

art by Andrew Macara,portrait of a sad woman, wearing a shirt with the text:\"No EGGS LEFT\"

- Model:  Stable Diffusion 3.5 Medium Turbo (SD3.5M Turbo).

- DPM++ 2M - Simple.

- 10 steps.

- LORAs: SD3.5M-Booster Type 1, SD3.5M-Booster Type 2, Samsung Galaxy S23 Ultra Photographic Style.

I started training loras for Flux, but recently I discovered that I could use all the datasets I used for Flux and use it again for SDXL and all things come out great because SDXL is so much lightier for training, as it is for inference, that I can put a lot more epoches and steps.

Now, when I go back to flux, it started to be a pain tô wait 10x more. For Flux I always used 16 or 8 epoches and for me it worked ok, but sometimes I fell flux do not learn details the way sdxl have been learning using 32 epoches, that is my current default for it (everything empirical).

So I have been wondering: would it worth training Flux for 32 epoches as well? Would it be a great improvement over 16 epoches?

I was going through the settings and I accidentally erased CFG Rescale and I can’t figure out how to get it back. I looked online but I’m only getting results on what it is. I know what it is it’s not telling me how to get it back. Any help is appreciated.

I've got an RTX 4090 and want to use it for AI video generation and training. I've got a concept in mind where I use StableDiffusion to generate images of characters, then I want to record video footage of myself in motion and speaking and overlay these generated characters over the top.

Is any open source software good enough for this purpose yet or will I need to buy something to get the results I want?

I use the draw things app and I use SDXL with a Pokemon trainer sprite LORA I found on civit. I can’t seem it figure out what’s going on but the line won’t go away

Hello, is there a guide pinned somewhere in this sub on how to make an ai model with photos and face swap videos?

Which would be better for trying out stable diffusion (can it also run flux/shuttle diffusion?) a $180 3060 or a $220 2080 ti? How much not having resize bar support effect the 2080 ti?

Hey all — throwback to a previous era. There used to be this *amazing* and comprehensive word document with tutorials for the Deforum Stable Diffusison notebook (local or colab). I can't seem to find it — anyone remember or know what I'm talking about by any chance? It used to have this gif on the opening page

https://preview.redd.it/jg5nwkgx739e1.png?width=120&format=png&auto=webp&s=9c7867925187fffc0cf090cd26b70783c05a68f7

Hi, I was using Last Ben Stable diffusion Git Hub (https://github.com/TheLastBen/fast-stable-diffusion).I have no knowledge of any software or code and not have a good laptop. Now this colab is showing error from lastweek (Screen shot attatched),all goes above my head. Any advice how to repair it or any other free colab would be appreciated. Thank you.

https://preview.redd.it/y0n6b4i2339e1.png?width=893&format=png&auto=webp&s=806e01e38f8b6f98f9ac28a5861341281b2eac42

https://preview.redd.it/w2jgb6j5339e1.png?width=1190&format=png&auto=webp&s=b05e0cb82ca64fb6ef881f08321260e36c821477

I would like to sign up for ModelsLab to use their text to video API and some others. They don't have a great reputation, judging by some of the online reviews, but there is also no other text to video service within my price point. Has anyone tried the $199 and $250 per month plan and if so how well do they scale? For my use case I'll probably need to generate a few thousand videos per month.